MagicAnimate: Temporally Consistent Human Image Animation using Diffusion Model
Remove background from an image
Falcons.ai Fine-Tuned Vision Transformer (ViT) for NSFW Image Classification
Implementation of Realistic Vision v5.1 with VAE
FLUX.1-Dev LoRA Explorer (DEPRECATED Please use: black-forest-labs/flux-dev-lora)
SDXL ControlNet - Canny
SDXL Inpainting by the HF Diffusers team
Robust face restoration algorithm for old photos/AI-generated faces
Juggernaut XL v9
Turn any image into a video
FLUX.1-Dev Multi LoRA Explorer
Segmind Stable Diffusion Model (SSD-1B) is a distilled 50% smaller version of SDXL, offering a 60% speedup while maintaining high-quality text-to-image generation capabilities
Hyper FLUX 8-step by ByteDance
FLUX.1-Schnell LoRA Explorer
CLIP Interrogator for SDXL optimizes text prompts to match a given image
Coqui XTTS-v2: Multilingual Text To Speech Voice Cloning
A multimodal LLM-based AI assistant, which is trained with alignment techniques. Qwen-VL-Chat supports more flexible interaction, such as multi-round question answering, and creative capabilities.
😊 Hotshot-XL is an AI text-to-GIF model trained to work alongside Stable Diffusion XL
SDXL v1.0 - A text-to-image generative AI model that creates beautiful images
snowflake-arctic-embed is a suite of text embedding models that focuses on creating high-quality retrieval models optimized for performance
Latent Consistency Model (LCM): SDXL, distills the original model into a version that requires fewer steps (4 to 8 instead of the original 25 to 50)
Monster Labs QrCode ControlNet on top of SD Realistic Vision v5.1
moondream2 is a small vision language model designed to run efficiently on edge devices
RealvisXL-v2.0 with LCM LoRA - requires fewer steps (4 to 8 instead of the original 40 to 50)
Implementation of SDXL RealVisXL_V2.0
Animate Your Personalized Text-to-Image Diffusion Models
Practical face restoration algorithm for *old photos* or *AI-generated faces* (for larger images)
DreamShaper is a general purpose SD model that aims at doing everything well, photos, art, anime, manga. It's designed to match Midjourney and DALL-E.
Real-ESRGAN Video Upscaler
Realistic Vision v5.0 Image 2 Image
A unique fusion that showcases exceptional prompt adherence and semantic understanding, it seems to be a step above base SDXL and a step closer to DALLE-3 in terms of prompt comprehension
CLIP Interrogator (for faster inference)
dreamshaper-xl-lightning is a Stable Diffusion model that has been fine-tuned on SDXL
AI-driven audio enhancement for your audio files, powered by Resemble AI
Phi-3-Mini-4K-Instruct is a 3.8B parameters, lightweight, state-of-the-art open model trained with the Phi-3 datasets
SDXL_Niji_Special Edition
Florence-2: Advancing a Unified Representation for a Variety of Vision Tasks
Realistic Vision V4.0
PixArt-Alpha 1024px is a transformer-based text-to-image diffusion system trained on text embeddings from T5
Ostris AI-Toolkit for Flux LoRA Training (DEPRECATED. Please use: ostris/flux-dev-lora-trainer)
Dreamshaper-7 img2img with LCM LoRA for faster inference
Florence-2: Advancing a Unified Representation for a Variety of Vision Tasks
Implementation of SDXL RealVisXL_V1.0
SDXL Image Blending
(Academic and Non-commercial use only) Pixel-Aware Stable Diffusion for Realistic Image Super-resolution and Personalized Stylization
BakLLaVA-1 is a Mistral 7B base augmented with the LLaVA 1.5 architecture
lmsys/vicuna-13b-v1.3
Realistic Vision v5.0 Inpainting
Mistral-7B-v0.1 fine tuned for chat with the Dolphin dataset (an open-source implementation of Microsoft's Orca)
Real-ESRGAN with optional face correction and adjustable upscale (for larger images)
Latest model in the Qwen family for chatting with video and image models
Realistic Vision v5.0 with VAE
Gemma2 2b by Google
The image prompt adapter is designed to enable a pretrained text-to-image diffusion model to generate SDXL images with an image prompt
SDXL ControlNet - Depth
(Research only) IP-Adapter-FaceID can generate various style images conditioned on a face with only text prompts
lmsys/vicuna-7b-v1.3
FalconAIs NSFW detection model, extended for videos
SDXL ControlNet - OpenPose
Meta's Llama 2 7b Chat - GPTQ
Segment Anything 2 (SAM2) by Meta - Automatic mask generation
Stylized Audio-Driven Single Image Talking Face Animation
sdxs-512-0.9 can generate high-resolution images in real-time based on prompt texts, trained using score distillation and feature matching
Meta's Llama 2 13b Chat - GPTQ
WizardCoder: Empowering Code Large Language Models with Evol-Instruct
ThinkDiffusionXL is a go-to model capable of amazing photorealism that's also versatile enough to generate high-quality images across a variety of styles and subjects without needing to be a prompting genius
This is wizard-vicuna-13b trained with a subset of the dataset - responses that contained alignment / moralizing were removed
Hyper FLUX 16-step by ByteDance
Ollama Llama 3.3 70B
Mistral-7B-v0.1 fine tuned for chat with the Dolphin dataset (an open-source implementation of Microsoft's Orca)
Image-to-video - SEINE: Short-to-Long Video Diffusion Model for Generative Transition and Prediction
Auto fuse a user's face onto the template image, with a similar appearance to the user
Segments an audio recording based on who is speaking (on A100)
(Research only) Moondream1 is a vision language model that performs on par with models twice its size
InterpAny-Clearer: Clearer anytime frame interpolation & Manipulated interpolation
Image to Image enhancer using DemoFusion
Open diffusion model for high-quality video generation
Orpheus 3B - high quality, emotive Text to Speech
SDXL lightning mult-controlnet, img2img & inpainting
DemoFusion: Democratising High-Resolution Image Generation With No 💰
Ollama Nemotron 70b
Implementation of SDXL RealVisXL_V2.0 img2img
Phi-3-Mini-128K-Instruct is a 3.8 billion-parameter, lightweight, state-of-the-art open model trained using the Phi-3 datasets
Stable Diffusion x4 upscaler model
In-Context LoRA with Image-to-Image and Inpainting to apply your logo to anything
Cog wrapper for Ollama llama3:70b
360 Panorama SDXL image with inpainted wrapping seam
Convert your videos to DensePose and use it with MagicAnimate
Projection module trained to add vision capabilties to Llama 3 using SigLIP
Realistic Vision V5 with OpenPose
Realistic Vision V3.0 with VAE
Fuyu-8B is a multi-modal text and image transformer trained by Adept AI
Qwen1.5 is the beta version of Qwen2, a transformer-based decoder-only language model pretrained on a large amount of data
Controlnet v1.1 - Tile Version
A combination of ip_adapter SDv1.5 and mediapipe-face to inpaint a face
SDXL using DeepCache
Segmind Stable Diffusion Model (SSD-1B) img2img
Phi-2 by Microsoft
Playground v2 is a diffusion-based text-to-image generative model trained from scratch. Try out all 3 models here
nomic-embed-text-v1 is 8192 context length text encoder that surpasses OpenAI text-embedding-ada-002 and text-embedding-3-small performance on short and long context tasks
