lucataco / sdxl-niji-se

Remove background from an image

Falcons.ai Fine-Tuned Vision Transformer (ViT) for NSFW Image Classification

Implementation of Realistic Vision v5.1 with VAE

FLUX.1-Dev LoRA Explorer (DEPRECATED Please use: black-forest-labs/flux-dev-lora)

SDXL ControlNet - Canny

SDXL Inpainting by the HF Diffusers team

Robust face restoration algorithm for old photos/AI-generated faces

Juggernaut XL v9

Turn any image into a video

FLUX.1-Dev Multi LoRA Explorer

Segmind Stable Diffusion Model (SSD-1B) is a distilled 50% smaller version of SDXL, offering a 60% speedup while maintaining high-quality text-to-image generation capabilities

Hyper FLUX 8-step by ByteDance

FLUX.1-Schnell LoRA Explorer

CLIP Interrogator for SDXL optimizes text prompts to match a given image

Coqui XTTS-v2: Multilingual Text To Speech Voice Cloning

A multimodal LLM-based AI assistant, which is trained with alignment techniques. Qwen-VL-Chat supports more flexible interaction, such as multi-round question answering, and creative capabilities.

😊 Hotshot-XL is an AI text-to-GIF model trained to work alongside Stable Diffusion XL

SDXL v1.0 - A text-to-image generative AI model that creates beautiful images

snowflake-arctic-embed is a suite of text embedding models that focuses on creating high-quality retrieval models optimized for performance

Latent Consistency Model (LCM): SDXL, distills the original model into a version that requires fewer steps (4 to 8 instead of the original 25 to 50)

Monster Labs QrCode ControlNet on top of SD Realistic Vision v5.1

moondream2 is a small vision language model designed to run efficiently on edge devices

RealvisXL-v2.0 with LCM LoRA - requires fewer steps (4 to 8 instead of the original 40 to 50)

Animate Your Personalized Text-to-Image Diffusion Models

Implementation of SDXL RealVisXL_V2.0

Animate Your Personalized Text-to-Image Diffusion Models

Practical face restoration algorithm for *old photos* or *AI-generated faces* (for larger images)

DreamShaper is a general purpose SD model that aims at doing everything well, photos, art, anime, manga. It's designed to match Midjourney and DALL-E.

Real-ESRGAN Video Upscaler

Realistic Vision v5.0 Image 2 Image

A unique fusion that showcases exceptional prompt adherence and semantic understanding, it seems to be a step above base SDXL and a step closer to DALLE-3 in terms of prompt comprehension

CLIP Interrogator (for faster inference)

dreamshaper-xl-lightning is a Stable Diffusion model that has been fine-tuned on SDXL

AI-driven audio enhancement for your audio files, powered by Resemble AI

Phi-3-Mini-4K-Instruct is a 3.8B parameters, lightweight, state-of-the-art open model trained with the Phi-3 datasets

Florence-2: Advancing a Unified Representation for a Variety of Vision Tasks

Realistic Vision V4.0

PixArt-Alpha 1024px is a transformer-based text-to-image diffusion system trained on text embeddings from T5

MagicAnimate: Temporally Consistent Human Image Animation using Diffusion Model

Ostris AI-Toolkit for Flux LoRA Training (DEPRECATED. Please use: ostris/flux-dev-lora-trainer)

Dreamshaper-7 img2img with LCM LoRA for faster inference

Florence-2: Advancing a Unified Representation for a Variety of Vision Tasks

Implementation of SDXL RealVisXL_V1.0

SDXL Image Blending

(Academic and Non-commercial use only) Pixel-Aware Stable Diffusion for Realistic Image Super-resolution and Personalized Stylization

BakLLaVA-1 is a Mistral 7B base augmented with the LLaVA 1.5 architecture

lmsys/vicuna-13b-v1.3

Realistic Vision v5.0 Inpainting

Mistral-7B-v0.1 fine tuned for chat with the Dolphin dataset (an open-source implementation of Microsoft's Orca)

Real-ESRGAN with optional face correction and adjustable upscale (for larger images)

Latest model in the Qwen family for chatting with video and image models

Realistic Vision v5.0 with VAE

Gemma2 2b by Google

The image prompt adapter is designed to enable a pretrained text-to-image diffusion model to generate SDXL images with an image prompt

SDXL ControlNet - Depth

(Research only) IP-Adapter-FaceID can generate various style images conditioned on a face with only text prompts

lmsys/vicuna-7b-v1.3

FalconAIs NSFW detection model, extended for videos

SDXL ControlNet - OpenPose

Meta's Llama 2 7b Chat - GPTQ

Segment Anything 2 (SAM2) by Meta - Automatic mask generation

Stylized Audio-Driven Single Image Talking Face Animation

sdxs-512-0.9 can generate high-resolution images in real-time based on prompt texts, trained using score distillation and feature matching

Meta's Llama 2 13b Chat - GPTQ

WizardCoder: Empowering Code Large Language Models with Evol-Instruct

ThinkDiffusionXL is a go-to model capable of amazing photorealism that's also versatile enough to generate high-quality images across a variety of styles and subjects without needing to be a prompting genius

This is wizard-vicuna-13b trained with a subset of the dataset - responses that contained alignment / moralizing were removed

Hyper FLUX 16-step by ByteDance

Ollama Llama 3.3 70B

Mistral-7B-v0.1 fine tuned for chat with the Dolphin dataset (an open-source implementation of Microsoft's Orca)

Image-to-video - SEINE: Short-to-Long Video Diffusion Model for Generative Transition and Prediction

Auto fuse a user's face onto the template image, with a similar appearance to the user

Segments an audio recording based on who is speaking (on A100)

(Research only) Moondream1 is a vision language model that performs on par with models twice its size

InterpAny-Clearer: Clearer anytime frame interpolation & Manipulated interpolation

Image to Image enhancer using DemoFusion

Open diffusion model for high-quality video generation

Orpheus 3B - high quality, emotive Text to Speech

SDXL lightning mult-controlnet, img2img & inpainting

DemoFusion: Democratising High-Resolution Image Generation With No 💰

Ollama Nemotron 70b

Implementation of SDXL RealVisXL_V2.0 img2img

Phi-3-Mini-128K-Instruct is a 3.8 billion-parameter, lightweight, state-of-the-art open model trained using the Phi-3 datasets

Stable Diffusion x4 upscaler model

In-Context LoRA with Image-to-Image and Inpainting to apply your logo to anything

Cog wrapper for Ollama llama3:70b

360 Panorama SDXL image with inpainted wrapping seam

Convert your videos to DensePose and use it with MagicAnimate

Projection module trained to add vision capabilties to Llama 3 using SigLIP

Realistic Vision V5 with OpenPose

Realistic Vision V3.0 with VAE

Fuyu-8B is a multi-modal text and image transformer trained by Adept AI

Qwen1.5 is the beta version of Qwen2, a transformer-based decoder-only language model pretrained on a large amount of data

Controlnet v1.1 - Tile Version

A combination of ip_adapter SDv1.5 and mediapipe-face to inpaint a face

SDXL using DeepCache

Segmind Stable Diffusion Model (SSD-1B) img2img

Phi-2 by Microsoft

Playground v2 is a diffusion-based text-to-image generative model trained from scratch. Try out all 3 models here

nomic-embed-text-v1 is 8192 context length text encoder that surpasses OpenAI text-embedding-ada-002 and text-embedding-3-small performance on short and long context tasks